Playground v2.5: Three Insights towards Enhancing Aesthetic Quality in
Text-to-Image Generation
- URL: http://arxiv.org/abs/2402.17245v1
- Date: Tue, 27 Feb 2024 06:31:52 GMT
- Title: Playground v2.5: Three Insights towards Enhancing Aesthetic Quality in
Text-to-Image Generation
- Authors: Daiqing Li, Aleks Kamko, Ehsan Akhgari, Ali Sabet, Linmiao Xu, Suhail
Doshi
- Abstract summary: We focus on enhancing color and contrast, improving generation across multiple aspect ratios, and improving human-centric fine details.
Our model is open-source, and we hope the development of Playground v2.5 provides valuable guidelines for researchers aiming to elevate the aesthetic quality of diffusion-based image generation models.
- Score: 3.976813869450304
- License: http://creativecommons.org/licenses/by/4.0/
- Abstract: In this work, we share three insights for achieving state-of-the-art
aesthetic quality in text-to-image generative models. We focus on three
critical aspects for model improvement: enhancing color and contrast, improving
generation across multiple aspect ratios, and improving human-centric fine
details. First, we delve into the significance of the noise schedule in
training a diffusion model, demonstrating its profound impact on realism and
visual fidelity. Second, we address the challenge of accommodating various
aspect ratios in image generation, emphasizing the importance of preparing a
balanced bucketed dataset. Lastly, we investigate the crucial role of aligning
model outputs with human preferences, ensuring that generated images resonate
with human perceptual expectations. Through extensive analysis and experiments,
Playground v2.5 demonstrates state-of-the-art performance in terms of aesthetic
quality under various conditions and aspect ratios, outperforming both
widely-used open-source models like SDXL and Playground v2, and closed-source
commercial systems such as DALLE 3 and Midjourney v5.2. Our model is
open-source, and we hope the development of Playground v2.5 provides valuable
guidelines for researchers aiming to elevate the aesthetic quality of
diffusion-based image generation models.
Related papers
- Illustrious: an Open Advanced Illustration Model [7.428509329724737]
We develop a text-to-image anime image generative model, called Illustrious, to achieve high resolution, dynamic color range images, and high restoration ability.
We focus on three critical approaches for model improvement. First, we delve into the significance of the batch size and dropout control, which enables faster learning of controllable token based concept activations.
Second, we increase the training resolution of images, affecting the accurate depiction of character anatomy in much higher resolution, extending its generation capability over 20MP with proper methods.
arXiv Detail & Related papers (2024-09-30T04:59:12Z) - FreeEnhance: Tuning-Free Image Enhancement via Content-Consistent Noising-and-Denoising Process [120.91393949012014]
FreeEnhance is a framework for content-consistent image enhancement using off-the-shelf image diffusion models.
In the noising stage, FreeEnhance is devised to add lighter noise to the region with higher frequency to preserve the high-frequent patterns in the original image.
In the denoising stage, we present three target properties as constraints to regularize the predicted noise, enhancing images with high acutance and high visual quality.
arXiv Detail & Related papers (2024-09-11T17:58:50Z) - Multimodal Large Language Model is a Human-Aligned Annotator for Text-to-Image Generation [87.50120181861362]
VisionPrefer is a high-quality and fine-grained preference dataset that captures multiple preference aspects.
We train a reward model VP-Score over VisionPrefer to guide the training of text-to-image generative models and the preference prediction accuracy of VP-Score is comparable to human annotators.
arXiv Detail & Related papers (2024-04-23T14:53:15Z) - YaART: Yet Another ART Rendering Technology [119.09155882164573]
This study introduces YaART, a novel production-grade text-to-image cascaded diffusion model aligned to human preferences.
We analyze how these choices affect both the efficiency of the training process and the quality of the generated images.
We demonstrate that models trained on smaller datasets of higher-quality images can successfully compete with those trained on larger datasets.
arXiv Detail & Related papers (2024-04-08T16:51:19Z) - Large-scale Reinforcement Learning for Diffusion Models [30.164571425479824]
Text-to-image diffusion models are susceptible to implicit biases that arise from web-scale text-image training pairs.
We present an effective scalable algorithm to improve diffusion models using Reinforcement Learning (RL)
We show how our approach substantially outperforms existing methods for aligning diffusion models with human preferences.
arXiv Detail & Related papers (2024-01-20T08:10:43Z) - Morphable Diffusion: 3D-Consistent Diffusion for Single-image Avatar Creation [14.064983137553353]
We aim to enhance the quality and functionality of generative diffusion models for the task of creating controllable, photorealistic human avatars.
We achieve this by integrating a 3D morphable model into the state-of-the-art multi-view-consistent diffusion approach.
Our proposed framework is the first diffusion model to enable the creation of fully 3D-consistent, animatable, and photorealistic human avatars.
arXiv Detail & Related papers (2024-01-09T18:59:04Z) - Kandinsky: an Improved Text-to-Image Synthesis with Image Prior and
Latent Diffusion [50.59261592343479]
We present Kandinsky1, a novel exploration of latent diffusion architecture.
The proposed model is trained separately to map text embeddings to image embeddings of CLIP.
We also deployed a user-friendly demo system that supports diverse generative modes such as text-to-image generation, image fusion, text and image fusion, image variations generation, and text-guided inpainting/outpainting.
arXiv Detail & Related papers (2023-10-05T12:29:41Z) - Breathing New Life into 3D Assets with Generative Repainting [74.80184575267106]
Diffusion-based text-to-image models ignited immense attention from the vision community, artists, and content creators.
Recent works have proposed various pipelines powered by the entanglement of diffusion models and neural fields.
We explore the power of pretrained 2D diffusion models and standard 3D neural radiance fields as independent, standalone tools.
Our pipeline accepts any legacy renderable geometry, such as textured or untextured meshes, and orchestrates the interaction between 2D generative refinement and 3D consistency enforcement tools.
arXiv Detail & Related papers (2023-09-15T16:34:51Z) - StyleAvatar3D: Leveraging Image-Text Diffusion Models for High-Fidelity
3D Avatar Generation [103.88928334431786]
We present a novel method for generating high-quality, stylized 3D avatars.
We use pre-trained image-text diffusion models for data generation and a Generative Adversarial Network (GAN)-based 3D generation network for training.
Our approach demonstrates superior performance over current state-of-the-art methods in terms of visual quality and diversity of the produced avatars.
arXiv Detail & Related papers (2023-05-30T13:09:21Z) - Quality-agnostic Image Captioning to Safely Assist People with Vision
Impairment [11.864465182761945]
We show how data augmentation techniques for generating synthetic noise can address data sparsity in this domain.
Second, we enhance the robustness of the model by expanding a state-of-the-art model to a dual network architecture.
Third, we evaluate the prediction reliability using confidence calibration on images with different difficulty/noise levels.
arXiv Detail & Related papers (2023-04-28T04:32:28Z)
This list is automatically generated from the titles and abstracts of the papers in this site.
This site does not guarantee the quality of this site (including all information) and is not responsible for any consequences.